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Feb 22, 2024 · The Stable Diffusion 3 suite of models currently ranges from 800M to 8B parameters. This approach aims to align with our core values and democratize access, providing users with a variety of options for scalability and quality to best meet their creative needs. Stable Diffusion 3 combines a diffusion transformer architecture and flow matching.
March 24, 2023. Stable UnCLIP 2.1. New stable diffusion finetune (Stable unCLIP 2.1, Hugging Face) at 768x768 resolution, based on SD2.1-768.This model allows for image variations and mixing operations as described in Hierarchical Text-Conditional Image Generation with CLIP Latents, and, thanks to its modularity, can be combined with other models such as KARLO.
Stable Diffusion Online. Stable Diffusion is a latent text-to-image diffusion model capable of generating photo-realistic images given any text input, cultivates autonomous freedom to produce incredible imagery, empowers billions of people to create stunning art within seconds.
Stable Diffusion is a deep learning, text-to-image model released in 2022 based on diffusion techniques. The generative artificial intelligence technology is the premier product of Stability AI and is considered to be a part of the ongoing artificial intelligence boom.. It is primarily used to generate detailed images conditioned on text descriptions, though it can also be applied to other ...
Stable Diffusion 3 Medium is the latest and most advanced text-to-image AI model in our Stable Diffusion 3 series, comprising two billion parameters. It excels in photorealism, processes complex prompts, and generates clear text. The weights are available under a community license.
Stable Diffusion is a Latent Diffusion model developed by researchers from the Machine Vision and Learning group at LMU Munich, a.k.a CompVis. Model checkpoints were publicly released at the end of August 2022 by a collaboration of Stability AI, CompVis, and Runway with support from EleutherAI and LAION.
Stable Diffusion 3 Medium is the latest and most advanced text-to-image AI model in our Stable Diffusion 3 series, comprising two billion parameters. It excels in photorealism, processes complex prompts, and generates clear text. The weights are available under a community license.
Note: Stable Diffusion v1 is a general text-to-image diffusion model and therefore mirrors biases and (mis-)conceptions that are present in its training data. Details on the training procedure and data, as well as the intended use of the model can be found in the corresponding model card.
Stable Diffusion generates image representation, a vector that numerically summarizes a high-resolution image depicted in the text prompt. This is done by refining a randomly initialized noise over multiple timesteps to gradually improve the image quality and adherence to the prompt.
Stable Diffusion is a latent diffusion model, a kind of deep generative neural network developed by the CompVis group at LMU Munich. The model has been released by a collaboration of Stability AI, CompVis LMU, and Runway with support from EleutherAI and LAION.
